1
0
mirror of https://github.com/huggingface/diffusers.git synced 2026-01-27 17:22:53 +03:00
Linoy Tsaban e30d3bf544 [LoRA] add LoRA support to HiDream and fine-tuning script (#11281)
* initial commit

* initial commit

* initial commit

* initial commit

* initial commit

* initial commit

* Update examples/dreambooth/train_dreambooth_lora_hidream.py

Co-authored-by: Bagheera <59658056+bghira@users.noreply.github.com>

* move prompt embeds, pooled embeds outside

* Update examples/dreambooth/train_dreambooth_lora_hidream.py

Co-authored-by: hlky <hlky@hlky.ac>

* Update examples/dreambooth/train_dreambooth_lora_hidream.py

Co-authored-by: hlky <hlky@hlky.ac>

* fix import

* fix import and tokenizer 4, text encoder 4 loading

* te

* prompt embeds

* fix naming

* shapes

* initial commit to add HiDreamImageLoraLoaderMixin

* fix init

* add tests

* loader

* fix model input

* add code example to readme

* fix default max length of text encoders

* prints

* nullify training cond in unpatchify for temp fix to incompatible shaping of transformer output during training

* smol fix

* unpatchify

* unpatchify

* fix validation

* flip pred and loss

* fix shift!!!

* revert unpatchify changes (for now)

* smol fix

* Apply style fixes

* workaround moe training

* workaround moe training

* remove prints

* to reduce some memory, keep vae in `weight_dtype` same as we have for flux (as it's the same vae)
bbd0c161b5/examples/dreambooth/train_dreambooth_lora_flux.py (L1207)

* refactor to align with HiDream refactor

* refactor to align with HiDream refactor

* refactor to align with HiDream refactor

* add support for cpu offloading of text encoders

* Apply style fixes

* adjust lr and rank for train example

* fix copies

* Apply style fixes

* update README

* update README

* update README

* fix license

* keep prompt2,3,4 as None in validation

* remove reverse ode comment

* Update examples/dreambooth/train_dreambooth_lora_hidream.py

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>

* Update examples/dreambooth/train_dreambooth_lora_hidream.py

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>

* vae offload change

* fix text encoder offloading

* Apply style fixes

* cleaner to_kwargs

* fix module name in copied from

* add requirements

* fix offloading

* fix offloading

* fix offloading

* update transformers version in reqs

* try AutoTokenizer

* try AutoTokenizer

* Apply style fixes

* empty commit

* Delete tests/lora/test_lora_layers_hidream.py

* change tokenizer_4 to load with AutoTokenizer as well

* make text_encoder_four and tokenizer_four configurable

* save model card

* save model card

* revert T5

* fix test

* remove non diffusers lumina2 conversion

---------

Co-authored-by: Bagheera <59658056+bghira@users.noreply.github.com>
Co-authored-by: hlky <hlky@hlky.ac>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: github-actions[bot] <github-actions[bot]@users.noreply.github.com>
2025-04-22 11:44:02 +03:00
2025-04-08 21:17:03 +05:30
2025-04-18 10:27:50 -10:00
2022-09-19 21:58:41 +02:00
2022-05-30 18:21:15 +02:00
2024-04-02 20:49:43 +05:30
2022-10-12 13:22:51 +02:00
2024-12-04 09:04:31 -08:00
2025-04-15 06:50:08 -10:00



GitHub GitHub release GitHub release Contributor Covenant X account

πŸ€— Diffusers is the go-to library for state-of-the-art pretrained diffusion models for generating images, audio, and even 3D structures of molecules. Whether you're looking for a simple inference solution or training your own diffusion models, πŸ€— Diffusers is a modular toolbox that supports both. Our library is designed with a focus on usability over performance, simple over easy, and customizability over abstractions.

πŸ€— Diffusers offers three core components:

  • State-of-the-art diffusion pipelines that can be run in inference with just a few lines of code.
  • Interchangeable noise schedulers for different diffusion speeds and output quality.
  • Pretrained models that can be used as building blocks, and combined with schedulers, for creating your own end-to-end diffusion systems.

Installation

We recommend installing πŸ€— Diffusers in a virtual environment from PyPI or Conda. For more details about installing PyTorch and Flax, please refer to their official documentation.

PyTorch

With pip (official package):

pip install --upgrade diffusers[torch]

With conda (maintained by the community):

conda install -c conda-forge diffusers

Flax

With pip (official package):

pip install --upgrade diffusers[flax]

Apple Silicon (M1/M2) support

Please refer to the How to use Stable Diffusion in Apple Silicon guide.

Quickstart

Generating outputs is super easy with πŸ€— Diffusers. To generate an image from text, use the from_pretrained method to load any pretrained diffusion model (browse the Hub for 30,000+ checkpoints):

from diffusers import DiffusionPipeline
import torch

pipeline = DiffusionPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", torch_dtype=torch.float16)
pipeline.to("cuda")
pipeline("An image of a squirrel in Picasso style").images[0]

You can also dig into the models and schedulers toolbox to build your own diffusion system:

from diffusers import DDPMScheduler, UNet2DModel
from PIL import Image
import torch

scheduler = DDPMScheduler.from_pretrained("google/ddpm-cat-256")
model = UNet2DModel.from_pretrained("google/ddpm-cat-256").to("cuda")
scheduler.set_timesteps(50)

sample_size = model.config.sample_size
noise = torch.randn((1, 3, sample_size, sample_size), device="cuda")
input = noise

for t in scheduler.timesteps:
    with torch.no_grad():
        noisy_residual = model(input, t).sample
        prev_noisy_sample = scheduler.step(noisy_residual, t, input).prev_sample
        input = prev_noisy_sample

image = (input / 2 + 0.5).clamp(0, 1)
image = image.cpu().permute(0, 2, 3, 1).numpy()[0]
image = Image.fromarray((image * 255).round().astype("uint8"))
image

Check out the Quickstart to launch your diffusion journey today!

How to navigate the documentation

Documentation What can I learn?
Tutorial A basic crash course for learning how to use the library's most important features like using models and schedulers to build your own diffusion system, and training your own diffusion model.
Loading Guides for how to load and configure all the components (pipelines, models, and schedulers) of the library, as well as how to use different schedulers.
Pipelines for inference Guides for how to use pipelines for different inference tasks, batched generation, controlling generated outputs and randomness, and how to contribute a pipeline to the library.
Optimization Guides for how to optimize your diffusion model to run faster and consume less memory.
Training Guides for how to train a diffusion model for different tasks with different training techniques.

Contribution

We ❀️ contributions from the open-source community! If you want to contribute to this library, please check out our Contribution guide. You can look out for issues you'd like to tackle to contribute to the library.

Also, say πŸ‘‹ in our public Discord channel Join us on Discord. We discuss the hottest trends about diffusion models, help each other with contributions, personal projects or just hang out β˜•.

Task Pipeline πŸ€— Hub
Unconditional Image Generation DDPM google/ddpm-ema-church-256
Text-to-Image Stable Diffusion Text-to-Image stable-diffusion-v1-5/stable-diffusion-v1-5
Text-to-Image unCLIP kakaobrain/karlo-v1-alpha
Text-to-Image DeepFloyd IF DeepFloyd/IF-I-XL-v1.0
Text-to-Image Kandinsky kandinsky-community/kandinsky-2-2-decoder
Text-guided Image-to-Image ControlNet lllyasviel/sd-controlnet-canny
Text-guided Image-to-Image InstructPix2Pix timbrooks/instruct-pix2pix
Text-guided Image-to-Image Stable Diffusion Image-to-Image stable-diffusion-v1-5/stable-diffusion-v1-5
Text-guided Image Inpainting Stable Diffusion Inpainting runwayml/stable-diffusion-inpainting
Image Variation Stable Diffusion Image Variation lambdalabs/sd-image-variations-diffusers
Super Resolution Stable Diffusion Upscale stabilityai/stable-diffusion-x4-upscaler
Super Resolution Stable Diffusion Latent Upscale stabilityai/sd-x2-latent-upscaler

Thank you for using us ❀️.

Credits

This library concretizes previous work by many different authors and would not have been possible without their great research and implementations. We'd like to thank, in particular, the following implementations which have helped us in our development and without which the API could not have been as polished today:

  • @CompVis' latent diffusion models library, available here
  • @hojonathanho original DDPM implementation, available here as well as the extremely useful translation into PyTorch by @pesser, available here
  • @ermongroup's DDIM implementation, available here
  • @yang-song's Score-VE and Score-VP implementations, available here

We also want to thank @heejkoo for the very helpful overview of papers, code and resources on diffusion models, available here as well as @crowsonkb and @rromb for useful discussions and insights.

Citation

@misc{von-platen-etal-2022-diffusers,
  author = {Patrick von Platen and Suraj Patil and Anton Lozhkov and Pedro Cuenca and Nathan Lambert and Kashif Rasul and Mishig Davaadorj and Dhruv Nair and Sayak Paul and William Berman and Yiyi Xu and Steven Liu and Thomas Wolf},
  title = {Diffusers: State-of-the-art diffusion models},
  year = {2022},
  publisher = {GitHub},
  journal = {GitHub repository},
  howpublished = {\url{https://github.com/huggingface/diffusers}}
}
Description
πŸ€— Diffusers: соврСмСнныС Π΄ΠΈΡ„Ρ„ΡƒΠ·ΠΈΠΎΠ½Π½Ρ‹Π΅ ΠΌΠΎΠ΄Π΅Π»ΠΈ для Π³Π΅Π½Π΅Ρ€Π°Ρ†ΠΈΠΈ ΠΈΠ·ΠΎΠ±Ρ€Π°ΠΆΠ΅Π½ΠΈΠΉ, Π²ΠΈΠ΄Π΅ΠΎ ΠΈ Π°ΡƒΠ΄ΠΈΠΎ Π² PyTorch.
Readme Apache-2.0 Cite this repository 754 MiB
Languages
Python 100%