* first add a script for DC-AE; * DC-AE init * replace triton with custom implementation * 1. rename file and remove un-used codes; * no longer rely on omegaconf and dataclass * replace custom activation with diffuers activation * remove dc_ae attention in attention_processor.py * iinherit from ModelMixin * inherit from ConfigMixin * dc-ae reduce to one file * update downsample and upsample * clean code * support DecoderOutput * remove get_same_padding and val2tuple * remove autocast and some assert * update ResBlock * remove contents within super().__init__ * Update src/diffusers/models/autoencoders/dc_ae.py Co-authored-by: YiYi Xu <yixu310@gmail.com> * remove opsequential * update other blocks to support the removal of build_norm * remove build encoder/decoder project in/out * remove inheritance of RMSNorm2d from LayerNorm * remove reset_parameters for RMSNorm2d Co-authored-by: YiYi Xu <yixu310@gmail.com> * remove device and dtype in RMSNorm2d __init__ Co-authored-by: YiYi Xu <yixu310@gmail.com> * Update src/diffusers/models/autoencoders/dc_ae.py Co-authored-by: YiYi Xu <yixu310@gmail.com> * Update src/diffusers/models/autoencoders/dc_ae.py Co-authored-by: YiYi Xu <yixu310@gmail.com> * Update src/diffusers/models/autoencoders/dc_ae.py Co-authored-by: YiYi Xu <yixu310@gmail.com> * remove op_list & build_block * remove build_stage_main * change file name to autoencoder_dc * move LiteMLA to attention.py * align with other vae decode output; * add DC-AE into init files; * update * make quality && make style; * quick push before dgx disappears again * update * make style * update * update * fix * refactor * refactor * refactor * update * possibly change to nn.Linear * refactor * make fix-copies * replace vae with ae * replace get_block_from_block_type to get_block * replace downsample_block_type from Conv to conv for consistency * add scaling factors * incorporate changes for all checkpoints * make style * move mla to attention processor file; split qkv conv to linears * refactor * add tests * from original file loader * add docs * add standard autoencoder methods * combine attention processor * fix tests * update * minor fix * minor fix * minor fix & in/out shortcut rename * minor fix * make style * fix paper link * update docs * update single file loading * make style * remove single file loading support; todo for DN6 * Apply suggestions from code review Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com> * add abstract --------- Co-authored-by: Junyu Chen <chenjydl2003@gmail.com> Co-authored-by: YiYi Xu <yixu310@gmail.com> Co-authored-by: chenjy2003 <70215701+chenjy2003@users.noreply.github.com> Co-authored-by: Aryan <aryan@huggingface.co> Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
π€ Diffusers is the go-to library for state-of-the-art pretrained diffusion models for generating images, audio, and even 3D structures of molecules. Whether you're looking for a simple inference solution or training your own diffusion models, π€ Diffusers is a modular toolbox that supports both. Our library is designed with a focus on usability over performance, simple over easy, and customizability over abstractions.
π€ Diffusers offers three core components:
- State-of-the-art diffusion pipelines that can be run in inference with just a few lines of code.
- Interchangeable noise schedulers for different diffusion speeds and output quality.
- Pretrained models that can be used as building blocks, and combined with schedulers, for creating your own end-to-end diffusion systems.
Installation
We recommend installing π€ Diffusers in a virtual environment from PyPI or Conda. For more details about installing PyTorch and Flax, please refer to their official documentation.
PyTorch
With pip (official package):
pip install --upgrade diffusers[torch]
With conda (maintained by the community):
conda install -c conda-forge diffusers
Flax
With pip (official package):
pip install --upgrade diffusers[flax]
Apple Silicon (M1/M2) support
Please refer to the How to use Stable Diffusion in Apple Silicon guide.
Quickstart
Generating outputs is super easy with π€ Diffusers. To generate an image from text, use the from_pretrained method to load any pretrained diffusion model (browse the Hub for 30,000+ checkpoints):
from diffusers import DiffusionPipeline
import torch
pipeline = DiffusionPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", torch_dtype=torch.float16)
pipeline.to("cuda")
pipeline("An image of a squirrel in Picasso style").images[0]
You can also dig into the models and schedulers toolbox to build your own diffusion system:
from diffusers import DDPMScheduler, UNet2DModel
from PIL import Image
import torch
scheduler = DDPMScheduler.from_pretrained("google/ddpm-cat-256")
model = UNet2DModel.from_pretrained("google/ddpm-cat-256").to("cuda")
scheduler.set_timesteps(50)
sample_size = model.config.sample_size
noise = torch.randn((1, 3, sample_size, sample_size), device="cuda")
input = noise
for t in scheduler.timesteps:
with torch.no_grad():
noisy_residual = model(input, t).sample
prev_noisy_sample = scheduler.step(noisy_residual, t, input).prev_sample
input = prev_noisy_sample
image = (input / 2 + 0.5).clamp(0, 1)
image = image.cpu().permute(0, 2, 3, 1).numpy()[0]
image = Image.fromarray((image * 255).round().astype("uint8"))
image
Check out the Quickstart to launch your diffusion journey today!
How to navigate the documentation
| Documentation | What can I learn? |
|---|---|
| Tutorial | A basic crash course for learning how to use the library's most important features like using models and schedulers to build your own diffusion system, and training your own diffusion model. |
| Loading | Guides for how to load and configure all the components (pipelines, models, and schedulers) of the library, as well as how to use different schedulers. |
| Pipelines for inference | Guides for how to use pipelines for different inference tasks, batched generation, controlling generated outputs and randomness, and how to contribute a pipeline to the library. |
| Optimization | Guides for how to optimize your diffusion model to run faster and consume less memory. |
| Training | Guides for how to train a diffusion model for different tasks with different training techniques. |
Contribution
We β€οΈ contributions from the open-source community! If you want to contribute to this library, please check out our Contribution guide. You can look out for issues you'd like to tackle to contribute to the library.
- See Good first issues for general opportunities to contribute
- See New model/pipeline to contribute exciting new diffusion models / diffusion pipelines
- See New scheduler
Also, say π in our public Discord channel . We discuss the hottest trends about diffusion models, help each other with contributions, personal projects or just hang out β.
Popular Tasks & Pipelines
| Task | Pipeline | π€ Hub |
|---|---|---|
| Unconditional Image Generation | DDPM | google/ddpm-ema-church-256 |
| Text-to-Image | Stable Diffusion Text-to-Image | stable-diffusion-v1-5/stable-diffusion-v1-5 |
| Text-to-Image | unCLIP | kakaobrain/karlo-v1-alpha |
| Text-to-Image | DeepFloyd IF | DeepFloyd/IF-I-XL-v1.0 |
| Text-to-Image | Kandinsky | kandinsky-community/kandinsky-2-2-decoder |
| Text-guided Image-to-Image | ControlNet | lllyasviel/sd-controlnet-canny |
| Text-guided Image-to-Image | InstructPix2Pix | timbrooks/instruct-pix2pix |
| Text-guided Image-to-Image | Stable Diffusion Image-to-Image | stable-diffusion-v1-5/stable-diffusion-v1-5 |
| Text-guided Image Inpainting | Stable Diffusion Inpainting | runwayml/stable-diffusion-inpainting |
| Image Variation | Stable Diffusion Image Variation | lambdalabs/sd-image-variations-diffusers |
| Super Resolution | Stable Diffusion Upscale | stabilityai/stable-diffusion-x4-upscaler |
| Super Resolution | Stable Diffusion Latent Upscale | stabilityai/sd-x2-latent-upscaler |
Popular libraries using 𧨠Diffusers
- https://github.com/microsoft/TaskMatrix
- https://github.com/invoke-ai/InvokeAI
- https://github.com/InstantID/InstantID
- https://github.com/apple/ml-stable-diffusion
- https://github.com/Sanster/lama-cleaner
- https://github.com/IDEA-Research/Grounded-Segment-Anything
- https://github.com/ashawkey/stable-dreamfusion
- https://github.com/deep-floyd/IF
- https://github.com/bentoml/BentoML
- https://github.com/bmaltais/kohya_ss
- +14,000 other amazing GitHub repositories πͺ
Thank you for using us β€οΈ.
Credits
This library concretizes previous work by many different authors and would not have been possible without their great research and implementations. We'd like to thank, in particular, the following implementations which have helped us in our development and without which the API could not have been as polished today:
- @CompVis' latent diffusion models library, available here
- @hojonathanho original DDPM implementation, available here as well as the extremely useful translation into PyTorch by @pesser, available here
- @ermongroup's DDIM implementation, available here
- @yang-song's Score-VE and Score-VP implementations, available here
We also want to thank @heejkoo for the very helpful overview of papers, code and resources on diffusion models, available here as well as @crowsonkb and @rromb for useful discussions and insights.
Citation
@misc{von-platen-etal-2022-diffusers,
author = {Patrick von Platen and Suraj Patil and Anton Lozhkov and Pedro Cuenca and Nathan Lambert and Kashif Rasul and Mishig Davaadorj and Dhruv Nair and Sayak Paul and William Berman and Yiyi Xu and Steven Liu and Thomas Wolf},
title = {Diffusers: State-of-the-art diffusion models},
year = {2022},
publisher = {GitHub},
journal = {GitHub repository},
howpublished = {\url{https://github.com/huggingface/diffusers}}
}
