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mirror of https://github.com/huggingface/diffusers.git synced 2026-01-27 17:22:53 +03:00
Sayak Paul 18fc40c169 [Feat] add tiny Autoencoder for (almost) instant decoding (#4384)
* add: model implementation of tiny autoencoder.

* add: inits.

* push the latest devs.

* add: conversion script and finish.

* add: scaling factor args.

* debugging

* fix denormalization.

* fix: positional argument.

* handle use_torch_2_0_or_xformers.

* handle post_quant_conv

* handle dtype

* fix: sdxl image processor for tiny ae.

* fix: sdxl image processor for tiny ae.

* unify upcasting logic.

* copied from madness.

* remove trailing whitespace.

* set is_tiny_vae = False

* address PR comments.

* change to AutoencoderTiny

* make act_fn an str throughout

* fix: apply_forward_hook decorator call

* get rid of the special is_tiny_vae flag.

* directly scale the output.

* fix dummies?

* fix: act_fn.

* get rid of the Clamp() layer.

* bring back copied from.

* movement of the blocks to appropriate modules.

* add: docstrings to AutoencoderTiny

* add: documentation.

* changes to the conversion script.

* add doc entry.

* settle tests.

* style

* add one slow test.

* fix

* fix 2

* fix 2

* fix: 4

* fix: 5

* finish integration tests

* Apply suggestions from code review

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* style

---------

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
2023-08-02 23:58:05 +05:30
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2022-09-19 21:58:41 +02:00
2023-03-14 20:16:30 -07:00
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2022-10-12 13:22:51 +02:00
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2023-07-26 23:06:18 +02:00



GitHub GitHub release Contributor Covenant

πŸ€— Diffusers is the go-to library for state-of-the-art pretrained diffusion models for generating images, audio, and even 3D structures of molecules. Whether you're looking for a simple inference solution or training your own diffusion models, πŸ€— Diffusers is a modular toolbox that supports both. Our library is designed with a focus on usability over performance, simple over easy, and customizability over abstractions.

πŸ€— Diffusers offers three core components:

  • State-of-the-art diffusion pipelines that can be run in inference with just a few lines of code.
  • Interchangeable noise schedulers for different diffusion speeds and output quality.
  • Pretrained models that can be used as building blocks, and combined with schedulers, for creating your own end-to-end diffusion systems.

Installation

We recommend installing πŸ€— Diffusers in a virtual environment from PyPi or Conda. For more details about installing PyTorch and Flax, please refer to their official documentation.

PyTorch

With pip (official package):

pip install --upgrade diffusers[torch]

With conda (maintained by the community):

conda install -c conda-forge diffusers

Flax

With pip (official package):

pip install --upgrade diffusers[flax]

Apple Silicon (M1/M2) support

Please refer to the How to use Stable Diffusion in Apple Silicon guide.

Quickstart

Generating outputs is super easy with πŸ€— Diffusers. To generate an image from text, use the from_pretrained method to load any pretrained diffusion model (browse the Hub for 4000+ checkpoints):

from diffusers import DiffusionPipeline
import torch

pipeline = DiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
pipeline.to("cuda")
pipeline("An image of a squirrel in Picasso style").images[0]

You can also dig into the models and schedulers toolbox to build your own diffusion system:

from diffusers import DDPMScheduler, UNet2DModel
from PIL import Image
import torch
import numpy as np

scheduler = DDPMScheduler.from_pretrained("google/ddpm-cat-256")
model = UNet2DModel.from_pretrained("google/ddpm-cat-256").to("cuda")
scheduler.set_timesteps(50)

sample_size = model.config.sample_size
noise = torch.randn((1, 3, sample_size, sample_size)).to("cuda")
input = noise

for t in scheduler.timesteps:
    with torch.no_grad():
        noisy_residual = model(input, t).sample
        prev_noisy_sample = scheduler.step(noisy_residual, t, input).prev_sample
        input = prev_noisy_sample

image = (input / 2 + 0.5).clamp(0, 1)
image = image.cpu().permute(0, 2, 3, 1).numpy()[0]
image = Image.fromarray((image * 255).round().astype("uint8"))
image

Check out the Quickstart to launch your diffusion journey today!

How to navigate the documentation

Documentation What can I learn?
Tutorial A basic crash course for learning how to use the library's most important features like using models and schedulers to build your own diffusion system, and training your own diffusion model.
Loading Guides for how to load and configure all the components (pipelines, models, and schedulers) of the library, as well as how to use different schedulers.
Pipelines for inference Guides for how to use pipelines for different inference tasks, batched generation, controlling generated outputs and randomness, and how to contribute a pipeline to the library.
Optimization Guides for how to optimize your diffusion model to run faster and consume less memory.
Training Guides for how to train a diffusion model for different tasks with different training techniques.

Contribution

We ❀️ contributions from the open-source community! If you want to contribute to this library, please check out our Contribution guide. You can look out for issues you'd like to tackle to contribute to the library.

Also, say πŸ‘‹ in our public Discord channel Join us on Discord. We discuss the hottest trends about diffusion models, help each other with contributions, personal projects or just hang out β˜•.

Task Pipeline πŸ€— Hub
Unconditional Image Generation DDPM google/ddpm-ema-church-256
Text-to-Image Stable Diffusion Text-to-Image runwayml/stable-diffusion-v1-5
Text-to-Image unclip kakaobrain/karlo-v1-alpha
Text-to-Image DeepFloyd IF DeepFloyd/IF-I-XL-v1.0
Text-to-Image Kandinsky kandinsky-community/kandinsky-2-2-decoder
Text-guided Image-to-Image Controlnet lllyasviel/sd-controlnet-canny
Text-guided Image-to-Image Instruct Pix2Pix timbrooks/instruct-pix2pix
Text-guided Image-to-Image Stable Diffusion Image-to-Image runwayml/stable-diffusion-v1-5
Text-guided Image Inpainting Stable Diffusion Inpaint runwayml/stable-diffusion-inpainting
Image Variation Stable Diffusion Image Variation lambdalabs/sd-image-variations-diffusers
Super Resolution Stable Diffusion Upscale stabilityai/stable-diffusion-x4-upscaler
Super Resolution Stable Diffusion Latent Upscale stabilityai/sd-x2-latent-upscaler

Thank you for using us ❀️

Credits

This library concretizes previous work by many different authors and would not have been possible without their great research and implementations. We'd like to thank, in particular, the following implementations which have helped us in our development and without which the API could not have been as polished today:

  • @CompVis' latent diffusion models library, available here
  • @hojonathanho original DDPM implementation, available here as well as the extremely useful translation into PyTorch by @pesser, available here
  • @ermongroup's DDIM implementation, available here
  • @yang-song's Score-VE and Score-VP implementations, available here

We also want to thank @heejkoo for the very helpful overview of papers, code and resources on diffusion models, available here as well as @crowsonkb and @rromb for useful discussions and insights.

Citation

@misc{von-platen-etal-2022-diffusers,
  author = {Patrick von Platen and Suraj Patil and Anton Lozhkov and Pedro Cuenca and Nathan Lambert and Kashif Rasul and Mishig Davaadorj and Thomas Wolf},
  title = {Diffusers: State-of-the-art diffusion models},
  year = {2022},
  publisher = {GitHub},
  journal = {GitHub repository},
  howpublished = {\url{https://github.com/huggingface/diffusers}}
}
Description
πŸ€— Diffusers: соврСмСнныС Π΄ΠΈΡ„Ρ„ΡƒΠ·ΠΈΠΎΠ½Π½Ρ‹Π΅ ΠΌΠΎΠ΄Π΅Π»ΠΈ для Π³Π΅Π½Π΅Ρ€Π°Ρ†ΠΈΠΈ ΠΈΠ·ΠΎΠ±Ρ€Π°ΠΆΠ΅Π½ΠΈΠΉ, Π²ΠΈΠ΄Π΅ΠΎ ΠΈ Π°ΡƒΠ΄ΠΈΠΎ Π² PyTorch.
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