mirror of
https://github.com/huggingface/diffusers.git
synced 2026-01-27 17:22:53 +03:00
[doc] add a tip about using SDXL refiner with hunyuan-dit and pixart (#8735)
* up * Apply suggestions from code review Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com> --------- Co-authored-by: Sayak Paul <spsayakpaul@gmail.com> Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
This commit is contained in:
@@ -34,6 +34,12 @@ Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.m
|
||||
|
||||
</Tip>
|
||||
|
||||
<Tip>
|
||||
|
||||
You can further improve generation quality by passing the generated image from [`HungyuanDiTPipeline`] to the [SDXL refiner](../../using-diffusers/sdxl#base-to-refiner-model) model.
|
||||
|
||||
</Tip>
|
||||
|
||||
## Optimization
|
||||
|
||||
You can optimize the pipeline's runtime and memory consumption with torch.compile and feed-forward chunking. To learn about other optimization methods, check out the [Speed up inference](../../optimization/fp16) and [Reduce memory usage](../../optimization/memory) guides.
|
||||
|
||||
@@ -37,6 +37,12 @@ Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers)
|
||||
|
||||
</Tip>
|
||||
|
||||
<Tip>
|
||||
|
||||
You can further improve generation quality by passing the generated image from [`PixArtSigmaPipeline`] to the [SDXL refiner](../../using-diffusers/sdxl#base-to-refiner-model) model.
|
||||
|
||||
</Tip>
|
||||
|
||||
## Inference with under 8GB GPU VRAM
|
||||
|
||||
Run the [`PixArtSigmaPipeline`] with under 8GB GPU VRAM by loading the text encoder in 8-bit precision. Let's walk through a full-fledged example.
|
||||
|
||||
@@ -285,6 +285,12 @@ refiner = DiffusionPipeline.from_pretrained(
|
||||
).to("cuda")
|
||||
```
|
||||
|
||||
<Tip warning={true}>
|
||||
|
||||
You can use SDXL refiner with a different base model. For example, you can use the [Hunyuan-DiT](../../api/pipelines/hunyuandit) or [PixArt-Sigma](../../api/pipelines/pixart_sigma) pipelines to generate images with better prompt adherence. Once you have generated an image, you can pass it to the SDXL refiner model to enhance final generation quality.
|
||||
|
||||
</Tip>
|
||||
|
||||
Generate an image from the base model, and set the model output to **latent** space:
|
||||
|
||||
```py
|
||||
|
||||
@@ -62,7 +62,7 @@ EXAMPLE_DOC_STRING = """
|
||||
>>> pipe = pipe.to(device)
|
||||
|
||||
>>> url = "https://raw.githubusercontent.com/CompVis/stable-diffusion/main/assets/stable-samples/img2img/sketch-mountains-input.jpg"
|
||||
>>> init_image = load_image(url).resize((512, 512))
|
||||
>>> init_image = load_image(url).resize((1024, 1024))
|
||||
|
||||
>>> prompt = "cat wizard, gandalf, lord of the rings, detailed, fantasy, cute, adorable, Pixar, Disney, 8k"
|
||||
|
||||
|
||||
Reference in New Issue
Block a user